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LearnSystem Design CapstonesDiffusion Models: Images & Text
👁️HardMultimodal Models

Diffusion Models: Images & Text

Design a governed image-generation service while learning DDPM noising, stochastic sampling, latent diffusion, classifier-free guidance, DiT backbones, and text diffusion.

49 min read
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Picture a designer asking for four previews of "oak dining table in soft light." The product promises a first preview below two seconds, but each variant still needs a sequence of denoiser evaluations before it can be decoded, checked, and published. At the required peak of 20 jobs per second, four variants and six denoising steps already mean 480 denoiser evaluations per second. If p95 slips, counting HTTP requests won't tell you whether the sampler, GPU, or release gate owns the delay.

That repeated repair loop is diffusion. Image models apply it to continuous pixels or latents; text models can apply a related idea to a discrete token canvas, refining masked positions instead of appending one token at a time.[1]Reference 1DiffusionGemmahttps://ai.google.dev/gemma/docs/diffusiongemma Start with one brightness value and compute its signal-to-noise mix by hand. Then follow that state through full images, latent pipelines, and discrete text diffusion.

Product brief: a governed design-asset generator

Design a service for product teams that creates synthetic campaign backgrounds and illustration previews. First decide where a slow or unsafe job can stop: before GPU work, inside the denoising loop, or after pixels are decoded. It must never present generated UI text, charts, or incident screenshots as verified evidence.

Requirements and API

  • POST /v1/image-jobs accepts prompt, negative_prompt, width, height, variant_count, route (preview or quality), seed, and an idempotency key. It returns job_id, accepted model and policy versions, and an estimated completion window.
  • GET /v1/image-jobs/{job_id} returns queued, running, review, failed, or complete, plus asset URLs, seed, sampler, step count, safety disposition, watermark result, and synthetic-content label.
  • Preview requests target p95 first preview below 2 seconds. Quality jobs may take longer but must finish within 30 seconds or return a visible failure.
  • Prompts and outputs pass policy checks. Logo use, realistic UI claims, and safety-sensitive imagery route to review. Tenant data and fine-tune adapters stay isolated.

Data flow and sizing

Store model, adapter, scheduler, policy, and seed versions with every asset so an operator can reproduce or explain it. Size the GPU loop in denoiser evaluations, not HTTP requests. A peak of 20 preview jobs/s with 4 variants and a 6-step distilled route needs 20 x 4 x 6 = 480 denoiser evaluations/s. That is the first capacity check, not a throughput claim for a particular GPU.

A smaller quality lane uses classifier-free guidance (CFG), which compares prompt-conditioned and unconditioned predictions at each step. At 2 jobs/s with 4 variants, 50 steps, and ordinary two-pass CFG, that lane needs 2 x 4 x 50 x 2 = 800 evaluations/s. Separate queues so offline quality work can't starve interactive previews.

Diagram showing Admission, GPU generation, Release, and POST /v1/image-jobs.
Admission, GPU generation, Release, and POST /v1/image-jobs.

Use the flow as an ownership map. Admission can reject a prompt before GPU spend. The GPU loop is where noise schedules, guidance, and step counts become latency. Release is where a pretty image can still be withheld. The rest of the chapter explains how to reason about each boundary.

Provenance belongs inside the asynchronous release path

A synthetic label in the API response is easy to lose after the file leaves the product. A durable release path embeds a detectable watermark in generated media, stores a separate provenance record, and verifies both before publication. Google DeepMind says SynthID embeds an imperceptible digital watermark when media is created. Its image and video watermark aims to remain detectable after modifications such as cropping, filters, and lossy compression.[2]Reference 2SynthIDhttps://deepmind.google/models/synthid/

That is a robustness target, not a cryptographic guarantee that survives every edit. Detection may be uncertain, and a missing watermark doesn't prove that the media is human-made. Treat the watermark as one provenance signal alongside content policy, metadata, and human review.

Diagram showing GPU worker decoded pixels, Output policy, Embed watermark, and Store private asset.
GPU worker decoded pixels, Output policy, Embed watermark, and Store private asset.

The job record should bind job_id, asset digest, model and adapter versions, seed, sampler, policy version, watermark method and version, detector version, detector result, and publication time. Keep the pre-watermark render private.

If resizing or transcoding happens later, run the detector on the original worker output and again on the final distributed derivative. The second check answers the only question that matters at publication: what will users actually receive?

The watermark stage must be idempotent. A worker retry either reuses the stored watermarked bytes or restarts from the stored generation configuration; it must not publish two assets with conflicting records under one job.

Measure watermark and detector latency separately from denoiser evaluations. Both happen after the GPU image loop, so a fast sampler won't hide a slow release path.

Why isn't synthetic=true response field enough provenance?

Answer

The field can disappear when the file is copied. Embed a detectable watermark, bind the asset digest to a durable job record, and verify the final distributed file before publishing. The absence of a detected watermark still doesn't prove human origin.

Failure recovery, rollout, and evaluation

Admission is idempotent, so a client retry returns the same job. A worker may rerun an unpublished job from its stored seed and versioned configuration.

A stochastic denoising diffusion probabilistic model (DDPM) retry must restore the random-number-generator state or restart the full sample. Resuming with fresh noise changes the asset. Timeout, out-of-memory, or failed safety checks move the job to failed or review rather than publishing a partial image.

Roll out a new model or sampler against a frozen prompt suite, then shadow it, canary 1% of eligible jobs, expand to 10%, and finally ramp by route. Keep the candidate and baseline on the same fixture set so a change in prompt mix doesn't masquerade as a model improvement.

Roll back on safety-policy regressions, memorization alerts, p95 latency breaches, or a material drop in human preference. Evaluation covers prompt adherence, human preference, diversity, text and layout fidelity, unsafe-output rate, memorization checks, p95 latency, failure rate, and cost per accepted asset.

A one-pixel forward process first

Before looking at a scheduler or a denoiser, answer one small question: at αˉ=0.5\bar{\alpha}=0.5αˉ=0.5, how much of a clean pixel remains? Fix one normalized pixel with brightness x0=1.0x_0 = 1.0x0​=1.0 and a known noise sample ϵ=−0.4\epsilon = -0.4ϵ=−0.4. At each noise level αˉ\bar{\alpha}αˉ, the closed-form mix is:

xt=αˉ x0+1−αˉ ϵx_t = \sqrt{\bar{\alpha}} \, x_0 + \sqrt{1 - \bar{\alpha}} \, \epsilonxt​=αˉ​x0​+1−αˉ​ϵ

αˉ\bar{\alpha}αˉSignal αˉ x0\sqrt{\bar{\alpha}}\,x_0αˉ​x0​Noise 1−αˉ ϵ\sqrt{1-\bar{\alpha}}\,\epsilon1−αˉ​ϵxtx_txt​SNR αˉ/(1−αˉ)\bar{\alpha}/(1-\bar{\alpha})αˉ/(1−αˉ)
0.90.949-0.1260.8229.0
0.50.707-0.2830.4241.0
0.0010.032-0.400-0.3680.001

At αˉ=0.5\bar{\alpha}=0.5αˉ=0.5, signal amplitude is 0.5 x0≈0.707 x0\sqrt{0.5}\,x_0 \approx 0.707\,x_00.5​x0​≈0.707x0​, not half brightness. Signal power is proportional to αˉ\bar{\alpha}αˉ, so half of the original power remains and SNR is 1.0.

At αˉ=0.001\bar{\alpha}=0.001αˉ=0.001, the pixel is almost pure noise. That is the forward story in miniature: as the noising timestep grows, the denoiser must recover structure from weaker signal.

If x0=1.0x_0=1.0x0​=1.0, ϵ=−0.4\epsilon=-0.4ϵ=−0.4, and αˉ=0.5\bar{\alpha}=0.5αˉ=0.5, what are xtx_txt​ and SNR?

Answer

Signal is 0.5≈0.707\sqrt{0.5}\approx0.7070.5​≈0.707, noise is 0.5⋅(−0.4)≈−0.283\sqrt{0.5}\cdot(-0.4)\approx-0.2830.5​⋅(−0.4)≈−0.283, so xt≈0.424x_t\approx0.424xt​≈0.424. SNR is 0.5/0.5=1.00.5/0.5=1.00.5/0.5=1.0.

adding-noise-gradually.py
1import json 2import math 3 4def forward_diffusion_value(x0: float, alpha_bar: float, epsilon: float) -> dict[str, float]: 5 signal = math.sqrt(alpha_bar) * x0 6 noise = math.sqrt(1 - alpha_bar) * epsilon 7 return { 8 "alpha_bar": alpha_bar, 9 "signal": round(signal, 3), 10 "noise": round(noise, 3), 11 "xt": round(signal + noise, 3), 12 "snr": round(alpha_bar / (1 - alpha_bar), 3), 13 } 14 15x0 = 1.0 16epsilon = -0.4 17steps = [ 18 forward_diffusion_value(x0, alpha_bar, epsilon) 19 for alpha_bar in [0.9, 0.5, 0.001] 20] 21 22print(json.dumps(steps, indent=2))
Output
1[ 2 { 3 "alpha_bar": 0.9, 4 "signal": 0.949, 5 "noise": -0.126, 6 "xt": 0.822, 7 "snr": 9.0 8 }, 9 { 10 "alpha_bar": 0.5, 11 "signal": 0.707, 12 "noise": -0.283, 13 "xt": 0.424, 14 "snr": 1.0 15 }, 16 { 17 "alpha_bar": 0.001, 18 "signal": 0.032, 19 "noise": -0.4, 20 "xt": -0.368, 21 "snr": 0.001 22 } 23]

Apply the same formula elementwise to every pixel or latent channel. Generation later reverses the path: it starts near pure noise and cleans structure back in step by step.[3]Reference 3Denoising Diffusion Probabilistic Models.https://arxiv.org/abs/2006.11239

Why the oak-table job isn't one generator pass

A one-shot mapping from a random vector to a 512x512 background sounds attractive. Generative Adversarial Networks (GANs) pair that generator with a discriminator. Their adversarial objective can be difficult to tune and can suffer mode collapse, while strong GAN implementations can still produce excellent samples. Diffusion changes the training objective and serving cost, not a guarantee of better campaign art.

The preview route doesn't ask the network for a clean image in one step. It asks for one small cleanup, like undoing one row of the table above. Repeating that cleanup turns a random field into a structured table and a room. The task at each step is simpler, the training objective is a supervised loss, and the prompt can steer every reverse step.

Classical DDPM training starts from an evidence lower bound (ELBO), a trainable lower bound on data log-likelihood. Its common simple objective is a noise-prediction regression loss derived from that bound.

Unlike GANs, this removes the adversarial minimax game. It doesn't by itself prove coverage, memorization safety, or output suitability.

Why did diffusion models become attractive compared with one-shot GAN-style generators?

Answer

They replace adversarial generator-discriminator optimization with supervised denoising targets and repeated conditioning opportunities. A production system still has to measure sample quality, diversity, memorization risk, and output safety on its intended use case.

Why this product pays for repeated denoiser steps

Count work before declaring a winner. GANs, Variational Autoencoders (VAEs), and diffusion models trade off sample quality, coverage, likelihood modeling, and inference cost. For campaign backgrounds, the costly column is inference work: one generator pass versus 6 preview steps or 50 quality steps. Treat the table as an architecture comparison, not a universal ranking:

FeatureGANsVAEsDiffusion Models
Training ObjectiveMinimax Game (Adversarial)ELBO (Evidence Lower Bound)Variational ELBO / MSE
Sample QualityCan be sharp; data and tuning dependentOften reconstruction-smoothedStrong results; model and sampler dependent
Coverage RiskMode collapse is an explicit riskCompression can remove detailMust evaluate diversity and memorization
Training BehaviorAdversarial instability riskStable reconstruction objectiveSupervised denoising objective
Inference WorkOne generator passOne decoder passRepeated denoiser evaluations unless accelerated
LikelihoodImplicitApproximateTractable lower bound

How the oak table becomes noise

Adding noise gradually

Ask what training must show the denoiser. It takes a clean campaign still, turns the table and room into isotropic noise, and asks the model to learn the reverse. The one-pixel table above is that story for one brightness value. The formal process extends it to full tensors and draws its principles from non-equilibrium thermodynamics.[4]Reference 4Deep Unsupervised Learning using Nonequilibrium Thermodynamics.https://arxiv.org/abs/1503.03585

Starting from a clean image x0x_0x0​, we add Gaussian noise over TTT steps. This fixed Markov chain means each new state depends only on the previous one, not on the full history. Structure gradually disappears until the data is indistinguishable from isotropic Gaussian noise, random static with the same variance on every pixel and color channel.[4]Reference 4Deep Unsupervised Learning using Nonequilibrium Thermodynamics.https://arxiv.org/abs/1503.03585

Forward diffusion visual with five image tiles from clean structure to terminal Gaussian noise, plus signal and noise bars showing signal fading as noise grows.
Forward diffusion gradually erases image structure while the noise term grows, and closed-form noising lets training jump straight to any timestep.

The transition kernel generalizes the one-pixel mix. Each step preserves a little signal and adds a little Gaussian noise:

q(xt∣xt−1)=N(xt;1−βt⋅xt−1,βtI)q(x_t | x_{t-1}) = \mathcal{N}(x_t; \sqrt{1-\beta_t} \cdot x_{t-1}, \beta_t \mathbf{I})q(xt​∣xt−1​)=N(xt​;1−βt​​⋅xt−1​,βt​I)

Here βt\beta_tβt​ is the noise schedule. The original DDPM experiments used a linear schedule from β1=10−4\beta_1 = 10^{-4}β1​=10−4 to βT=0.02\beta_T = 0.02βT​=0.02 over T=1000T = 1000T=1000 steps; newer systems may choose different schedules.[3]Reference 3Denoising Diffusion Probabilistic Models.https://arxiv.org/abs/2006.11239 Because sums of Gaussians remain Gaussian, the derivation can marginalize over intermediate steps. Let αt=1−βt\alpha_t = 1 - \beta_tαt​=1−βt​ and αˉt=∏s=1tαs\bar{\alpha}_t = \prod_{s=1}^{t} \alpha_sαˉt​=∏s=1t​αs​:

q(xt∣x0)=N(xt;αˉt⋅x0,(1−αˉt)I)q(x_t | x_0) = \mathcal{N}(x_t; \sqrt{\bar{\alpha}_t} \cdot x_0, (1 - \bar{\alpha}_t) \mathbf{I})q(xt​∣x0​)=N(xt​;αˉt​​⋅x0​,(1−αˉt​)I)

This is the formula you computed by hand. Sample xtx_txt​ at any timestep directly from x0x_0x0​ and one noise draw without simulating intermediate states.

At t=Tt = Tt=T, αˉT≈0\bar{\alpha}_T \approx 0αˉT​≈0, so xT≈N(0,I)x_T \approx \mathcal{N}(0, \mathbf{I})xT​≈N(0,I), pure isotropic Gaussian noise.

Why is the closed-form forward sample useful during training?

Answer

It lets the trainer pick any timestep independently, add the exact amount of Gaussian noise in one operation, and train batches across many noise levels in parallel. Without it, each training sample would have to simulate every previous forward step.

At a small forward timestep ttt, the input retains recognizable structure. At a large forward ttt near TTT, it retains very little original signal. The SNR curve αˉt/(1−αˉt)\bar{\alpha}_t / (1 - \bar{\alpha}_t)αˉt​/(1−αˉt​) describes that spectrum. Keep its direction separate from reverse sampling: forward time adds noise, while reverse time removes it.

Reverse process: learned denoising

DDPM training

Training has one job: approximate the reverse process pθ(xt−1∣xt)p_\theta(x_{t-1} | x_t)pθ​(xt−1​∣xt​), the step-by-step cleanup that generation will run later.[3]Reference 3Denoising Diffusion Probabilistic Models.https://arxiv.org/abs/2006.11239

Why not compute that reverse step exactly? The posterior q(xt−1∣xt)q(x_{t-1} \mid x_t)q(xt−1​∣xt​) marginalizes over every clean image x0x_0x0​ that could have produced xtx_txt​, and that marginalization makes it intractable.

The conditional posterior q(xt−1∣xt,x0)q(x_{t-1} \mid x_t, x_0)q(xt−1​∣xt​,x0​) is tractable and remains Gaussian. DDPM uses that closed form in its derivation, then trains a model to predict enough information from xtx_txt​ and ttt to approximate the reverse step without seeing x0x_0x0​ at inference time.[3]Reference 3Denoising Diffusion Probabilistic Models.https://arxiv.org/abs/2006.11239

Why is the reverse process learned instead of computed exactly?

Answer

Given only a noisy image xtx_txt​, there are many possible clean images that could have produced it, so q(xt−1∣xt)q(x_{t-1} \mid x_t)q(xt−1​∣xt​) is not directly available. Training gives the denoiser a learned estimate of the noise, clean image, or velocity needed by the scheduler to approximate the reverse step.

The model doesn't have to predict the clean image x0x_0x0​ directly. In standard DDPM (Denoising Diffusion Probabilistic Models), it predicts the noise component ϵ\epsilonϵ added at step ttt.

Train a neural network ϵθ(xt,t)\epsilon_\theta(x_t, t)ϵθ​(xt​,t) to recover that noise from a noisy input:

Lsimple=Et,x0,ϵ[∥ϵ−ϵθ(xt,t)∥2]\mathcal{L}_{\text{simple}} = \mathbb{E}_{t, x_0, \epsilon}[\|\epsilon - \epsilon_\theta(x_t, t)\|^2]Lsimple​=Et,x0​,ϵ​[∥ϵ−ϵθ​(xt​,t)∥2]

For each row, sample a timestep ttt, a clean image x0x_0x0​, and noise ϵ\epsilonϵ. Build xtx_txt​, ask the network to predict ϵ\epsilonϵ, and average the squared error over many rows.

In practice, the "simple" objective is an MSE loss on predicted noise. It comes from simplifying and reweighting the original ELBO, but the implementation is supervised regression from xtx_txt​ to ϵ\epsilonϵ.[3]Reference 3Denoising Diffusion Probabilistic Models.https://arxiv.org/abs/2006.11239 The small program below makes that target explicit:

ddpm-training.py
1import json 2 3def mse(true_values: list[float], predicted_values: list[float]) -> float: 4 squared_errors = [ 5 (true - predicted) ** 2 6 for true, predicted in zip(true_values, predicted_values, strict=True) 7 ] 8 return sum(squared_errors) / len(squared_errors) 9 10training_rows = [ 11 { 12 "timestep": 100, 13 "true_noise": [0.20, -0.10, 0.05], 14 "predicted_noise": [0.18, -0.08, 0.02], 15 }, 16 { 17 "timestep": 700, 18 "true_noise": [0.90, -0.40, 0.30], 19 "predicted_noise": [0.72, -0.35, 0.10], 20 }, 21] 22 23losses = [ 24 { 25 "timestep": row["timestep"], 26 "mse_loss": round(mse(row["true_noise"], row["predicted_noise"]), 4), 27 } 28 for row in training_rows 29] 30 31print(json.dumps(losses, indent=2))
Output
1[ 2 { 3 "timestep": 100, 4 "mse_loss": 0.0006 5 }, 6 { 7 "timestep": 700, 8 "mse_loss": 0.025 9 } 10]

The first row has a small error because its predicted noise is close to the injected noise. The second row is harder. Across training, the same network must handle every timestep, from nearly clean images to nearly pure static.

There is no adversarial game and no discriminator. The ground-truth target is available because training knows exactly which ϵ\epsilonϵ it injected.

What target does the standard DDPM training loop regress against?

Answer

It predicts the exact Gaussian noise ϵ\epsilonϵ that was added to the clean sample at timestep ttt. The loss is usually mean squared error between true noise and predicted noise.

Noise schedule and SNR

The schedule isn't just bookkeeping. Predict what the denoiser should see at a high-SNR versus a low-SNR timestep. In the closed-form forward equation, the signal coefficient is αˉt\sqrt{\bar{\alpha}_t}αˉt​​ and the noise coefficient is 1−αˉt\sqrt{1 - \bar{\alpha}_t}1−αˉt​​, so the effective signal-to-noise ratio at timestep ttt is:

SNR(t)=αˉt1−αˉt\mathrm{SNR}(t) = \frac{\bar{\alpha}_t}{1 - \bar{\alpha}_t}SNR(t)=1−αˉt​αˉt​​

High-SNR timesteps still contain recognizable structure. Low-SNR timesteps are close to pure noise.

If the schedule destroys signal too quickly, many late timesteps become almost uninformative and the model spends capacity on inputs with little semantic content left. Schedule design and loss weighting therefore decide where the model spends its learning budget along the SNR curve.

Forward-mix coefficients versus alpha-bar from 0 to 1. Signal sqrt(alpha-bar) rises toward 1 while noise sqrt(1-alpha-bar) falls toward 0. Markers at 0.001, 0.5, and 0.9 match the one-pixel table: 0.032/1.0, 0.707/0.707, and 0.949/0.316.
The x-axis is alpha-bar, not evenly spaced checkpoints. At 0.5 both amplitudes are 0.707 and SNR is 1.0, matching the one-pixel table.

What does the noise schedule control besides bookkeeping?

Answer

It controls the SNR curve across timesteps. That decides whether the model spends learning capacity on high-signal cleanup, low-signal structure hallucination, or a useful mix of both.

Prediction targets

Original DDPM predicts ϵ\epsilonϵ directly.[3]Reference 3Denoising Diffusion Probabilistic Models.https://arxiv.org/abs/2006.11239 That isn't the only valid parameterization. Some systems predict the clean sample x0x_0x0​; others predict a velocity-style target vvv, a linear mix such as v=αˉ ϵ−1−αˉ x0v = \sqrt{\bar{\alpha}}\,\epsilon - \sqrt{1-\bar{\alpha}}\,x_0v=αˉ​ϵ−1−αˉ​x0​ (progressive-distillation / flow-adjacent form).[5]Reference 5Progressive Distillation for Fast Sampling of Diffusion Models.https://arxiv.org/abs/2202.00512

These targets are mathematically convertible, but they change optimization behavior and how errors spread across timesteps. The scheduler still defines the reverse update; the prediction target defines what the network learns to output.

If a model predicts vvv instead of ϵ\epsilonϵ, did the diffusion pipeline stop being diffusion?

Answer

No. ϵ\epsilonϵ, x0x_0x0​, and vvv are different parameterizations used by the denoiser and scheduler. The pipeline still iteratively updates a noisy state toward a clean sample; the target changes optimization and error distribution, not the core reverse-process idea.

Sampling (reverse process)

To generate a new image, start from pure noise xT∼N(0,I)x_T \sim \mathcal{N}(0, \mathbf{I})xT​∼N(0,I) and iteratively denoise it. A DDPM reverse step combines a learned mean with fresh stochastic noise on non-final steps:

xt−1=1αt(xt−βt1−αˉtϵθ(xt,t))+σtz,z∼N(0,I)\begin{aligned} x_{t-1} &= \frac{1}{\sqrt{\alpha_t}} \left(x_t - \frac{\beta_t}{\sqrt{1-\bar{\alpha}_t}} \epsilon_\theta(x_t,t)\right) + \sigma_t z, \\ z &\sim \mathcal{N}(0,\mathbf{I}) \end{aligned}xt−1​z​=αt​​1​(xt​−1−αˉt​​βt​​ϵθ​(xt​,t))+σt​z,∼N(0,I)​

Set z=0z=0z=0 for the final step into x0x_0x0​. The original formulation allows a chosen reverse variance; the scalar teaching example uses σt=βt\sigma_t=\sqrt{\beta_t}σt​=βt​​.

Each step depends on the previous state, so the reverse process remains sequential even when denoiser calls are batched. A seeded random generator makes the stochastic trace reproducible:

sampling-reverse-process.py
1import json 2import math 3import random 4 5betas = [0.02, 0.03, 0.04] 6alphas = [1 - beta for beta in betas] 7alpha_bars: list[float] = [] 8running = 1.0 9for alpha in alphas: 10 running *= alpha 11 alpha_bars.append(running) 12 13def predicted_noise(x: float, step: int) -> float: 14 return 0.45 * x + 0.02 * step 15 16x = 1.25 17trace = [] 18rng = random.Random(7) 19for step in reversed(range(len(betas))): 20 beta_t = betas[step] 21 alpha_t = alphas[step] 22 alpha_bar_t = alpha_bars[step] 23 eps = predicted_noise(x, step) 24 mean = (x - (beta_t / math.sqrt(1 - alpha_bar_t)) * eps) / math.sqrt(alpha_t) 25 z = rng.gauss(0, 1) if step > 0 else 0.0 26 x = mean + math.sqrt(beta_t) * z 27 trace.append({ 28 "step": step, 29 "predicted_noise": round(eps, 3), 30 "reverse_mean": round(mean, 3), 31 "stochastic_z": round(z, 3), 32 "next_x": round(x, 3), 33 }) 34 35print(json.dumps(trace, indent=2))
Output
1[ 2 { 3 "step": 2, 4 "predicted_noise": 0.603, 5 "reverse_mean": 1.193, 6 "stochastic_z": -0.256, 7 "next_x": 1.141 8 }, 9 { 10 "step": 1, 11 "predicted_noise": 0.534, 12 "reverse_mean": 1.086, 13 "stochastic_z": 0.511, 14 "next_x": 1.174 15 }, 16 { 17 "step": 0, 18 "predicted_noise": 0.528, 19 "reverse_mean": 1.111, 20 "stochastic_z": 0.0, 21 "next_x": 1.111 22 } 23]

Read the trace as a dependency chain: each new next_x becomes the next step's input, while stochastic_z changes intermediate states. The final step sets it to zero.

The original DDPM sampler evaluates the denoiser across its configured reverse chain, commonly presented with T=1000T=1000T=1000 steps. That baseline is expensive compared with one-pass generation, which motivates sparse samplers and distilled models. Before taking those shortcuts, we need a denoiser that can predict the noise and a way to inject the prompt into every reverse step.

Model architectures: U-Net and diffusion transformers

The U-Net backbone

A denoiser must recover both layout and small spatial details. For a UI illustration scene, it needs to place panels and visual anchors coherently while preserving edges, textures, and lighting cues. Ask what happens if it builds layout but loses local edges: the scene may be plausible while every label and border is soft.

A U-Net (U-shaped convolutional network with skip connections) addresses that failure with two paths. Its down path builds broader context, its up path restores spatial resolution, and skip connections carry local cues around the deepest bottleneck.

In text-conditioned image models, timestep embeddings modulate residual blocks, and cross-attention is inserted at several resolutions so prompt tokens can steer both coarse layout and fine detail.[6]Reference 6High-Resolution Image Synthesis with Latent Diffusion Models.https://arxiv.org/abs/2112.10752

U-Net denoiser path showing down blocks building context, a mid block, up blocks restoring detail, skip connections for local edges, and timestep plus text conditioning at down, mid, and up blocks.
U-Net builds context on the way down, restores detail on the way up, and keeps local cues through skip paths.

Read these as representative conditioning sites, not every block. Stable Diffusion-style U-Nets use cross-attention in multiple down, middle, and up blocks rather than only at the bottleneck.[6]Reference 6High-Resolution Image Synthesis with Latent Diffusion Models.https://arxiv.org/abs/2112.10752

Why do U-Nets fit denoising tasks well?

Answer

They combine coarse scene context from the down path with spatial detail from skip connections on the up path. That matches denoising: early layers reason about layout while later layers restore local structure and texture.

The transformer shift

A Diffusion Transformer (DiT) changes the denoiser's representation. It patchifies the input into tokens and processes them with transformer blocks instead of a convolutional U-Net.

In the DiT experiments, increasing model compute improved Fréchet Inception Distance (FID), an image-quality metric where lower is better, for latent diffusion settings.[7]Reference 7Scalable Diffusion Models with Transformers.https://arxiv.org/abs/2212.09748 That supports a scaling path. It doesn't promise that a DiT is cheaper or better for every resolution, latency budget, or dataset.

Stable Diffusion 3's MM-DiT (multimodal diffusion transformer) is one published example. Its text and image representations participate in a joint attention operation while retaining separate modality-specific weights, rather than feeding text only through one-directional cross-attention.[8]Reference 8Scaling Rectified Flow Transformers for High-Resolution Image Synthesishttps://arxiv.org/abs/2403.03206

The SD3 authors report prompt-following, typography, and scaling results for that recipe. Those measured results shouldn't be generalized to every transformer-based image generator.

What changes when a denoiser moves from U-Net to DiT?

Answer

The noisy latent is patchified into tokens and processed by transformer blocks instead of convolutional encoder-decoder blocks. DiT reports favorable scaling results, but attention cost still depends on token count. Stable Diffusion 3 is a published example of a multimodal DiT with joint text-image attention; other model families need their own documentation and evaluation.

Timestep conditioning

Both backbones need to know how much noise remains. U-Net-style models usually encode the timestep ttt with a sinusoidal embedding, project it with an MLP, and use the result to modulate residual blocks through FiLM (Feature-wise Linear Modulation) or AdaGN (Adaptive Group Normalization).

DiT carries the same timestep signal, but typically injects it through adaptive LayerNorm inside transformer blocks instead of GroupNorm-based residual blocks.[7]Reference 7Scalable Diffusion Models with Transformers.https://arxiv.org/abs/2212.09748 A framework-free scale-and-shift sketch makes the shared idea visible:

timestep-conditioning.py
1import json 2 3def modulate(features: list[float], scale: float, shift: float) -> list[float]: 4 return [round(value * (1 + scale) + shift, 3) for value in features] 5 6feature_map = [0.20, -0.10, 0.50] 7timestep_controls = { 8 "reverse_start_high_noise": {"scale": 0.50, "shift": -0.05}, 9 "reverse_end_low_noise": {"scale": 0.05, "shift": 0.02}, 10} 11 12outputs = { 13 name: modulate(feature_map, **control) 14 for name, control in timestep_controls.items() 15} 16 17print(json.dumps(outputs, indent=2))
Output
1{ 2 "reverse_start_high_noise": [ 3 0.25, 4 -0.2, 5 0.7 6 ], 7 "reverse_end_low_noise": [ 8 0.23, 9 -0.085, 10 0.545 11 ] 12}

This modulation acts like a global controller. Reverse sampling starts at large ttt with high noise, where the network establishes broad structure. It ends near t=0t=0t=0 with low noise, where the network refines texture and detail.

In the forward noising direction, those same endpoints are called large or late ttt and small or early ttt, respectively. Name the direction when discussing "early" so the step is unambiguous.

Why must the denoiser know the timestep?

Answer

The same noisy tensor can require different behavior depending on how much signal remains. Timestep conditioning tells the model whether it should recover broad structure from heavy noise or refine local details near the end.

Flow matching and rectified flow

DDPM frames generation as reversing a stochastic noising chain. Flow matching reframes it as learning a velocity field that transports samples along a path between a noise distribution and the data distribution.[9]Reference 9Flow Matching for Generative Modelinghttps://arxiv.org/abs/2210.02747

Rectified flow picks a simple path: a straight line interpolating noise x1∼N(0,I)x_1 \sim \mathcal{N}(0, \mathbf{I})x1​∼N(0,I) and data x0x_0x0​ as xτ=(1−τ)x0+τx1x_\tau = (1 - \tau) x_0 + \tau x_1xτ​=(1−τ)x0​+τx1​ for τ∈[0,1]\tau \in [0, 1]τ∈[0,1].[10]Reference 10Flow Straight and Fast: Learning to Generate and Transfer Data with Rectified Flowhttps://arxiv.org/abs/2209.03003 The network learns the constant velocity x1−x0x_1 - x_0x1​−x0​ along that line, another regression target much like ϵ\epsilonϵ or vvv:

LFM=Eτ,x0,x1[ ∥ vθ(xτ,τ)−(x1−x0) ∥2 ]\mathcal{L}_{\text{FM}} = \mathbb{E}_{\tau, x_0, x_1}\big[\,\|\,v_\theta(x_\tau, \tau) - (x_1 - x_0)\,\|^2\,\big]LFM​=Eτ,x0​,x1​​[∥vθ​(xτ​,τ)−(x1​−x0​)∥2]

Straight-line paths aim to make numerical integration effective with fewer solver steps, but quality at any step count remains a measurement, not a guarantee. Stable Diffusion 3 reports rectified flow together with MM-DiT and a shifted timestep-sampling strategy.[8]Reference 8Scaling Rectified Flow Transformers for High-Resolution Image Synthesishttps://arxiv.org/abs/2403.03206

In a design review, connect flow matching to the same deployment concerns as diffusion: conditioning, latent representation, and sampler cost. Keep the training objective and scheduler equations distinct.

How does rectified flow differ from a standard DDPM noise schedule?

Answer

DDPM reverses a stochastic noising chain and predicts the added noise. Rectified flow defines a straight-line path between noise and data and trains the network to predict velocity along that line. Straight paths are intended to integrate effectively with fewer solver steps; Stable Diffusion 3 reports this objective with its MM-DiT backbone, but delivered quality and latency still need measurement.

Text diffusion and DiffusionGemma

Image diffusion corrupts continuous pixel or latent values with Gaussian noise. Text is different: a sentence is a sequence of discrete token IDs. A text diffuser therefore uses a discrete corruption process, often replacing tokens with masks or sampling from categorical transitions, then trains a model to reconstruct original tokens from noisy context.[11]Reference 11Structured Denoising Diffusion Models in Discrete State-Spaceshttps://arxiv.org/abs/2107.03006

Autoregressive language models generate left to right: p(xi∣x<i)p(x_i \mid x_{<i})p(xi​∣x<i​). Once token 12 is emitted, token 4 usually won't be revised.

A diffusion language model starts from a noisy or masked sequence and refines many positions over several rounds. Each round can use context from both sides of a missing token, so the model can repair an earlier word after seeing later words.

Google's DiffusionGemma is a current open-weights example. Google's docs describe it as based on the Gemma 4 26B-A4B mixture-of-experts family. The model card reports 25.2B total parameters and 3.8B active, with 8 of 128 experts plus one shared expert.[1]Reference 1DiffusionGemmahttps://ai.google.dev/gemma/docs/diffusiongemma[12]Reference 12DiffusionGemma Model Cardhttps://ai.google.dev/gemma/docs/diffusiongemma/model_card

It accepts text, image, and video inputs and generates text with an encoder-decoder stack. An autoregressive encoder prefills the prompt into a KV cache. A decoder then applies bidirectional attention over a 256-token generation canvas and iteratively denoises that block. When a canvas is finished, the encoder appends it to the cache before the next canvas starts. This is block-autoregressive discrete diffusion, not a single pass over an entire document or ordinary next-token decoding.[12]Reference 12DiffusionGemma Model Cardhttps://ai.google.dev/gemma/docs/diffusiongemma/model_card

Recommended sampling in the public docs uses a 48-step upper bound, adaptive early stopping that typically finishes in about 12 to 16 steps, and an entropy-bound sampler that commits only tokens the model is relatively sure about.[1]Reference 1DiffusionGemmahttps://ai.google.dev/gemma/docs/diffusiongemma

Google also notes that the speed advantage is strongest at low-to-medium batch sizes on a single accelerator. High-QPS cloud batching can shrink that gap.

Predict the visible difference before looking at the figure: left-to-right decoding commits a prefix, while a diffusion canvas can change several positions before it commits.

Autoregressive decoding appends one token at a time while discrete text diffusion starts from masked token slots and refines many positions over repeated passes before committing the final sequence.
Autoregressive decoding commits left to right, while text diffusion refines a whole token canvas over several passes.

The important systems difference is dependency shape. Autoregressive decoding has a hard sequential dependency between adjacent generated tokens. Text diffusion still needs multiple model passes, but each pass can update many uncertain positions in one block.

That makes it attractive for high-throughput short or medium outputs, infilling, and editing workflows where bidirectional context matters. It also introduces new engineering questions: how long should the canvas be, how many refinement rounds are enough, and how should the UI handle text that can change across the whole block before it's final?

The toy example below isn't a language model. It shows the state shape: a masked canvas gets refined in rounds, and later rounds can fill earlier positions after seeing surrounding context.

text-diffusion-refinement.py
1import json 2 3target = ["write", "unit", "tests", "first"] 4canvas = ["<mask>"] * len(target) 5update_plan = [ 6 {"round": 1, "positions": [0, 3]}, 7 {"round": 2, "positions": [1]}, 8 {"round": 3, "positions": [2]}, 9] 10 11trace = [] 12for step in update_plan: 13 for position in step["positions"]: 14 canvas[position] = target[position] 15 trace.append({ 16 "round": step["round"], 17 "canvas": " ".join(canvas), 18 "remaining_masks": canvas.count("<mask>"), 19 }) 20 21print(json.dumps(trace, indent=2))
Output
1[ 2 { 3 "round": 1, 4 "canvas": "write <mask> <mask> first", 5 "remaining_masks": 2 6 }, 7 { 8 "round": 2, 9 "canvas": "write unit <mask> first", 10 "remaining_masks": 1 11 }, 12 { 13 "round": 3, 14 "canvas": "write unit tests first", 15 "remaining_masks": 0 16 } 17]

Text diffusion shouldn't be described as "Stable Diffusion but with words." Its forward corruption is discrete, its scheduler decides which token positions remain uncertain, and its output length often has to be planned up front. The shared idea is denoising; data type, conditioning, and serving constraints differ.

What changes when diffusion moves from images to text?

Answer

The noisy state becomes a discrete token sequence rather than a continuous image or latent tensor. Generation can refine many masked positions over repeated passes, which changes latency and editing tradeoffs, but it also requires canvas-length planning and different corruption/scheduling rules.

Classifier-free guidance (CFG)

How CFG steers generation without a classifier

How does the model know to draw an "oak dining table" and not a "red running shoe"? Early conditional diffusion systems used classifier guidance. They trained a separate image classifier on noisy inputs, then added the classifier's gradient into the sampling step to push the image toward a target label. That worked, but it required training and maintaining an extra classifier.

CFG (Classifier-Free Guidance) removes that extra classifier. One denoiser learns both conditional and unconditional predictions, then inference extrapolates from the unconditional direction toward the text-conditioned one.[13]Reference 13Classifier-Free Diffusion Guidance.https://arxiv.org/abs/2207.12598

What problem does classifier-free guidance solve?

Answer

It strengthens prompt adherence without maintaining a separate noisy-image classifier. One denoiser learns conditional and unconditional predictions, then inference extrapolates from the unconditional direction toward the text-conditioned direction.

During training, randomly drop the conditioning, such as a text prompt, with probability puncondp_{uncond}puncond​. Ho and Salimans compared 10%, 20%, and 50% on 64x64 ImageNet and found 10% and 20% worked about equally well.[13]Reference 13Classifier-Free Diffusion Guidance.https://arxiv.org/abs/2207.12598

In real text-to-image systems, the unconditional branch is usually an empty-prompt embedding or learned null token, not a literal all-zero vector. The demo uses a higher drop rate so a four-row batch shows both branches:

cfg-training-step.py
1import json 2import random 3 4def cfg_conditioning_batch(prompts: list[str], drop_prob: float, seed: int = 7) -> list[dict[str, str]]: 5 rng = random.Random(seed) 6 rows = [] 7 for prompt in prompts: 8 dropped = rng.random() < drop_prob 9 rows.append({ 10 "original_prompt": prompt, 11 "conditioning_used": "<null>" if dropped else prompt, 12 "branch": "unconditional" if dropped else "conditional", 13 }) 14 return rows 15 16batch = cfg_conditioning_batch( 17 [ 18 "oak dining table in soft light", 19 "oak dining table in natural light", 20 "mobile dashboard empty state", 21 "annotated latency chart hero image", 22 ], 23 drop_prob=0.35, 24) 25 26print(json.dumps(batch, indent=2))
Output
1[ 2 { 3 "original_prompt": "oak dining table in soft light", 4 "conditioning_used": "<null>", 5 "branch": "unconditional" 6 }, 7 { 8 "original_prompt": "oak dining table in natural light", 9 "conditioning_used": "<null>", 10 "branch": "unconditional" 11 }, 12 { 13 "original_prompt": "mobile dashboard empty state", 14 "conditioning_used": "mobile dashboard empty state", 15 "branch": "conditional" 16 }, 17 { 18 "original_prompt": "annotated latency chart hero image", 19 "conditioning_used": "<null>", 20 "branch": "unconditional" 21 } 22]

During sampling, combine the conditional and unconditional predictions. Common guidance scales above 111 extrapolate beyond the conditional prediction:

ϵ^=ϵθ(xt,t,∅)+w⋅(ϵθ(xt,t,c)−ϵθ(xt,t,∅))\hat{\epsilon} = \epsilon_\theta(x_t, t, \emptyset) + w \cdot (\epsilon_\theta(x_t, t, c) - \epsilon_\theta(x_t, t, \emptyset))ϵ^=ϵθ​(xt​,t,∅)+w⋅(ϵθ​(xt​,t,c)−ϵθ​(xt​,t,∅))

Here, www is the guidance scale. When w=1w=1w=1, the equation reduces to the conditional prediction ϵθ(xt,t,c)\epsilon_\theta(x_t, t, c)ϵθ​(xt​,t,c).

When w>1w>1w>1, the model amplifies the difference between conditional and unconditional predictions. It pushes the image away from generic features and toward the text prompt.

Watch the convention. This is the guidance-scale form used in Stable Diffusion and the diffusers library, where w=1w=1w=1 means no extra steering. The original paper writes the same thing as ϵ^=(1+w′)ϵθ(xt,t,c)−w′ϵθ(xt,t,∅)\hat{\epsilon} = (1 + w')\epsilon_\theta(x_t, t, c) - w' \epsilon_\theta(x_t, t, \emptyset)ϵ^=(1+w′)ϵθ​(xt​,t,c)−w′ϵθ​(xt​,t,∅), where w′=0w'=0w′=0 means no steering.[13]Reference 13Classifier-Free Diffusion Guidance.https://arxiv.org/abs/2207.12598 The two map exactly via w=1+w′w = 1 + w'w=1+w′, so a reported guidance scale of 7.57.57.5 here corresponds to w′=6.5w' = 6.5w′=6.5 in the paper's notation.

A concrete way to read the formula: if the unconditional prediction says "this patch looks like generic wood grain" and the conditional prediction says "this patch looks like oak grain," w=7.5w = 7.5w=7.5 pushes the result seven and a half times farther along that direction. The patch becomes more distinctly oak-like rather than generic wood.

The practical implementation below shows CFG extrapolation on short vectors. Real systems can batch unconditional and conditional inputs into one forward pass, then split the predictions before applying the same math.

cfg-guidance-formula.py
1import json 2 3def guided_noise( 4 unconditional: list[float], 5 conditional: list[float], 6 guidance_scale: float, 7) -> list[float]: 8 return [ 9 round(base + guidance_scale * (target - base), 3) 10 for base, target in zip(unconditional, conditional, strict=True) 11 ] 12 13noise_uncond = [0.20, -0.10, 0.05] 14noise_cond = [0.05, -0.35, 0.20] 15 16rows = [ 17 { 18 "guidance_scale": scale, 19 "guided_noise": guided_noise(noise_uncond, noise_cond, scale), 20 } 21 for scale in [0.0, 1.0, 7.5, 15.0] 22] 23 24print(json.dumps(rows, indent=2))
Output
1[ 2 { 3 "guidance_scale": 0.0, 4 "guided_noise": [ 5 0.2, 6 -0.1, 7 0.05 8 ] 9 }, 10 { 11 "guidance_scale": 1.0, 12 "guided_noise": [ 13 0.05, 14 -0.35, 15 0.2 16 ] 17 }, 18 { 19 "guidance_scale": 7.5, 20 "guided_noise": [ 21 -0.925, 22 -1.975, 23 1.175 24 ] 25 }, 26 { 27 "guidance_scale": 15.0, 28 "guided_noise": [ 29 -2.05, 30 -3.85, 31 2.3 32 ] 33 } 34]
Guidance Scale wwwEffect
0.0Unconditional sample. Ignores the prompt entirely.
1.0Pure conditional prediction. Good baseline, no extra CFG boost.
3.0 to 8.0Illustrative tuning range. Measure adherence, diversity, and artifacts for this model and route.
15.0+Over-guided. Higher artifact risk, oversaturation, and repeated textures.
Classifier-free guidance as vector extrapolation on the article noise example. Unconditional noise sits at (0.20, -0.10), conditional noise at (0.05, -0.35), and guidance scale 7.5 continues that difference to (-0.925, -1.975).
The first two noise components from the worked example. w=1 stops at the conditional point; w=7.5 keeps going along the same difference vector.

Why is guidance scale not a probability?

Answer

It's an extrapolation coefficient between unconditional and conditional noise predictions. Larger values push harder toward prompt features, but they can reduce diversity and introduce oversaturation, texture repetition, or anatomy artifacts.

Ordinary CFG still represents two denoiser predictions per step, even when a system batches them into one kernel launch. Distillation can remove one branch at runtime. That's why sampler choice, distillation, and consistency models show up next: they attack evaluation count, not HTTP request count.

Faster sampling with implicit diffusion

Denoising Diffusion Implicit Models (DDIM) reuse a diffusion model's training objective while sampling across a selected subset of timesteps. The practical question is whether fewer evaluations preserve the route's required detail and control. A deterministic setting can also avoid fresh noise at each reverse step.

Before reading the routes below, predict the tradeoff: a sparse schedule should lower latency, while a distilled or consistency model changes what the network learned.

Representative generation routes showing a 1000-evaluation DDPM reference chain, a 50-evaluation DDIM timestep subset, a 6-evaluation distilled student, and a 4-evaluation consistency model.
DDIM selects a sparse timestep subset, while distilled and consistency routes change what the model learns. Counts are representative route budgets, not fixed method limits. The 6-step distilled row matches the preview-lane sizing in the product brief.

Compare denoiser-evaluation budgets, not just labels. The original DDPM setup is commonly demonstrated with about 1000 reverse steps. A DDIM deployment may try a sparse schedule such as 20 or 50 steps, while distilled or consistency models target still fewer.

Whether a reduced schedule preserves acceptable quality depends on model, prompts, guidance, and evaluation slice.

DDIM constructs a non-Markovian forward family with the same per-timestep marginals and training objective used by DDPM. Its reverse trajectory still updates one selected state into the next selected state.

The practical difference is a shorter timestep subsequence, with one setting making that trajectory deterministic.[14]Reference 14Denoising Diffusion Implicit Models.https://arxiv.org/abs/2010.02502

DDIM lets the sampler skip steps. Instead of sampling t,t−1,t−2t, t-1, t-2t,t−1,t−2, it can step t,t−Δ,t−2Δt, t-\Delta, t-2\Deltat,t−Δ,t−2Δ. This can reduce denoiser evaluations substantially.

It doesn't eliminate the need to compare prompt adherence, detail preservation, and edit behavior against the slower baseline.

xt−1=αˉt−1⋅x^0⏟predicted x0+1−αˉt−1−σt2⋅ϵθ(xt,t)+σt⋅z\begin{aligned} x_{t-1} &= \sqrt{\bar{\alpha}_{t-1}} \cdot \underbrace{\hat{x}_0}_{\text{predicted } x_0} \\ &\quad + \sqrt{1 - \bar{\alpha}_{t-1} - \sigma_t^2} \cdot \epsilon_\theta(x_t, t) \\ &\quad + \sigma_t \cdot z \end{aligned}xt−1​​=αˉt−1​​⋅predicted x0​x^0​​​+1−αˉt−1​−σt2​​⋅ϵθ​(xt​,t)+σt​⋅z​

Here x^0\hat{x}_0x^0​ is the model's clean-image estimate, ϵθ(xt,t)\epsilon_\theta(x_t,t)ϵθ​(xt​,t) is the predicted noise direction, and σtz\sigma_t zσt​z injects randomness. That term vanishes for deterministic DDIM when σt=0\sigma_t = 0σt​=0.

When the network is trained for ϵ\epsilonϵ-prediction (as in the DDPM objective above), recover the clean-image estimate before the DDIM update:

x^0=xt−1−αˉt ϵθ(xt,t)αˉt\hat{x}_0 = \frac{x_t - \sqrt{1-\bar{\alpha}_t}\,\epsilon_\theta(x_t,t)}{\sqrt{\bar{\alpha}_t}}x^0​=αˉt​​xt​−1−αˉt​​ϵθ​(xt​,t)​

That identity bridges the training target ϵθ\epsilon_\thetaϵθ​ to the x^0\hat{x}_0x^0​ term in the sampler equation. Indices ttt and t−1t-1t−1 in the display are adjacent points on the selected schedule; they need not be unit steps when you skip.

Setting σt=0\sigma_t = 0σt​=0 makes the trajectory deterministic. By bypassing the strict Markovian requirement, DDIM decoupled generation from the exact number of forward steps used during training.

You can train on 1000 noise levels but sample on a sparse subset during deployment. Training and inference schedules don't need to be identical.

What is the practical systems insight behind DDIM?

Answer

The model can train over many noise levels but sample on a sparse subset at deployment. Skipping steps reduces latency, while deterministic updates make image inversion and reproducible edits easier.

Each step can require more than one denoiser evaluation. For ordinary CFG, a system needs conditional and unconditional predictions unless it batches or distills that work. Count evaluations before claiming an interactive route is affordable:

sampler-evaluation-budget.py
1import json 2 3routes = [ 4 {"name": "offline_ddim_cfg", "steps": 50, "predictions_per_step": 2}, 5 {"name": "interactive_distilled", "steps": 6, "predictions_per_step": 1}, 6] 7 8budget = [ 9 { 10 "route": route["name"], 11 "denoiser_evaluations": route["steps"] * route["predictions_per_step"], 12 "quality_gate_required": True, 13 } 14 for route in routes 15] 16 17print(json.dumps(budget, indent=2))
Output
1[ 2 { 3 "route": "offline_ddim_cfg", 4 "denoiser_evaluations": 100, 5 "quality_gate_required": true 6 }, 7 { 8 "route": "interactive_distilled", 9 "denoiser_evaluations": 6, 10 "quality_gate_required": true 11 } 12]

Latent diffusion

Architecture

Suppose the denoiser takes the same number of steps in two representations. Pixel-space diffusion processes high-resolution feature maps such as 512×512×3512 \times 512 \times 3512×512×3 at every step. Latent Diffusion Models (LDMs), like Stable Diffusion, move that loop into the latent space of a high-capacity autoencoder.

A variational autoencoder has two parts: an encoder that compresses an image into a compact latent representation and a decoder that reconstructs the image from it. This is learned lossy compression. It preserves enough perceptual structure for the training objective, while reconstruction errors remain part of final output quality.[6]Reference 6High-Resolution Image Synthesis with Latent Diffusion Models.https://arxiv.org/abs/2112.10752

This is the continuous-latent branch of the autoencoder family. A VQ-VAE uses a discrete codebook instead: each patch is snapped to a learned code ID, making the compressed image look more like a sequence of tokens.[15]Reference 15Neural Discrete Representation Learning.https://arxiv.org/abs/1711.00937 That design suits a Transformer trained to predict image tokens autoregressively or with masking.

Latent diffusion usually takes the continuous route. It adds noise to the compressed tensor, then trains a denoiser over that continuous state.

During training, the VAE encoder maps an image to a clean latent z0z_0z0​, and the model learns to denoise noisy versions of that latent.

During inference, there's no input image. Start from random latent noise zT∼N(0,I)z_T \sim \mathcal{N}(0, \mathbf{I})zT​∼N(0,I), run the reverse process in latent space, and decode once at the end.[6]Reference 6High-Resolution Image Synthesis with Latent Diffusion Models.https://arxiv.org/abs/2112.10752

Latent diffusion path where prompt conditioning is reused across denoising steps, noisy latent grids become structured, and one final VAE decode produces pixels.
Prompt conditioning is reused at every denoising step. Pixels appear only after one final VAE decode of the clean latent.
Latent diffusion compares training and inference in compressed latent space, showing one VAE encode, many latent denoising steps, and one final VAE decode back to pixels.
Training encodes a real image once, then denoises known noise. Inference starts from random latents and decodes pixels once.

Only the final denoised latent gets decoded. If the denoiser predicts noise or velocity, the scheduler uses that prediction to update zt→zt−1z_t \rightarrow z_{t-1}zt​→zt−1​. The decoder never consumes raw predicted noise.

Why latent space?

Moving diffusion from pixels to latents splits the workload into two phases. The autoencoder compresses and reconstructs pixels with measurable detail loss. The diffusion model performs repeated denoising on the smaller latent state.

A 64×64×464 \times 64 \times 464×64×4 latent tensor has 48×48\times48× fewer scalar values than a 512×512×3512 \times 512 \times 3512×512×3 image. That is the Stable Diffusion 1.x-style shape used in the worked count below, not a universal latent.

Later systems can change channel count or spatial compression; Stable Diffusion 3 reports a 16-channel autoencoder.[8]Reference 8Scaling Rectified Flow Transformers for High-Resolution Image Synthesishttps://arxiv.org/abs/2403.03206 Count the latent values the loop actually iterates on, then benchmark the resulting route.

Why does Stable Diffusion-style latent space change production economics?

Answer

The expensive iterative loop runs on a compressed latent tensor instead of full pixels. That cuts spatial work substantially, while the VAE decoder reconstructs pixels once at the end. Benchmark throughput and inspect reconstruction loss on the details the route must preserve.

The count gives three checks. A 64×64×464 \times 64 \times 464×64×4 latent tensor has 48×48\times48× fewer scalar values than a 512×512×3512 \times 512 \times 3512×512×3 image and 64×64\times64× fewer spatial locations. That smaller state reduces denoising work relative to pixel-space diffusion.[6]Reference 6High-Resolution Image Synthesis with Latent Diffusion Models.https://arxiv.org/abs/2112.10752

It also removes information. Evaluate reconstruction quality on the details the route cares about.

latent-work-state.py
1import json 2import math 3 4pixel_shape = (512, 512, 3) 5latent_shape = (64, 64, 4) 6pixel_values = math.prod(pixel_shape) 7latent_values = math.prod(latent_shape) 8 9print(json.dumps({ 10 "pixel_values": pixel_values, 11 "latent_values": latent_values, 12 "scalar_reduction": pixel_values // latent_values, 13 "spatial_reduction": (pixel_shape[0] * pixel_shape[1]) // (latent_shape[0] * latent_shape[1]), 14 "throughput_requires_benchmark": True, 15}, indent=2))
Output
1{ 2 "pixel_values": 786432, 3 "latent_values": 16384, 4 "scalar_reduction": 48, 5 "spatial_reduction": 64, 6 "throughput_requires_benchmark": true 7}

Text conditioning via cross-attention

Text-conditioning recipes are model-specific. The latent diffusion paper demonstrates cross-attention conditioning and uses pretrained encoders such as CLIP for text-to-image experiments.[6]Reference 6High-Resolution Image Synthesis with Latent Diffusion Models.https://arxiv.org/abs/2112.10752[16]Reference 16Learning Transferable Visual Models From Natural Language Supervision.https://arxiv.org/abs/2103.00020 Other image generators may use T5-family encoders, multiple text encoders, or different freeze/train choices.

The stable concept is that prompt embeddings condition repeated denoising updates. It isn't that every text encoder is frozen.

The CrossAttentionBlock takes flattened image or latent features as queries, and text embeddings as keys and values. A numerically stable softmax turns compatibility scores into weights. The runnable version shows one latent patch attending to prompt tokens:

text-conditioning-via-cross-attention.py
1import json 2import math 3 4def dot(left: list[float], right: list[float]) -> float: 5 return sum(a * b for a, b in zip(left, right, strict=True)) 6 7def softmax(values: list[float]) -> list[float]: 8 largest = max(values) 9 exp_values = [math.exp(value - largest) for value in values] 10 total = sum(exp_values) 11 return [value / total for value in exp_values] 12 13latent_query = [0.25, 0.10, 0.70] 14text_tokens = [ 15 {"token": "oak", "key": [0.30, 0.05, 0.65], "value": [0.70, 0.20]}, 16 {"token": "dining", "key": [0.10, 0.80, 0.10], "value": [0.20, 0.70]}, 17 {"token": "table", "key": [0.35, 0.10, 0.55], "value": [0.80, 0.10]}, 18] 19 20logits = [ 21 dot(latent_query, token["key"]) / math.sqrt(len(latent_query)) 22 for token in text_tokens 23] 24weights = softmax(logits) 25conditioned_patch = [ 26 sum(weight * token["value"][i] for weight, token in zip(weights, text_tokens, strict=True)) 27 for i in range(2) 28] 29 30print(json.dumps({ 31 "attention_weights": { 32 token["token"]: round(weight, 3) 33 for token, weight in zip(text_tokens, weights, strict=True) 34 }, 35 "conditioned_patch": [round(value, 3) for value in conditioned_patch], 36}, indent=2))
Output
1{ 2 "attention_weights": { 3 "oak": 0.359, 4 "dining": 0.292, 5 "table": 0.349 6 }, 7 "conditioned_patch": [ 8 0.589, 9 0.311 10 ] 11}

In text-to-image cross-attention, what supplies queries, keys, and values?

Answer

Image or latent features usually supply the queries. Text-token embeddings supply keys and values, so each spatial location can attend to prompt concepts relevant to layout, objects, style, and attributes.

Control and specialization

Once the base route works, teams usually add targeted controls around it instead of retraining the whole backbone.

  • LoRA (Low-Rank Adaptation): Full fine-tuning is expensive. LoRA[17]Reference 17LoRA: Low-Rank Adaptation of Large Language Models.https://arxiv.org/abs/2106.09685 inserts small trainable matrices into selected layers so you can specialize a base model with a much smaller parameter update.
  • Inpainting and outpainting: The model conditions on known pixels and a mask, which lets it fill holes, replace objects, or extend the canvas while preserving local context.
  • ControlNet / adapters: ControlNet (a neural network architecture for adding conditional control to diffusion models)[18]Reference 18Adding Conditional Control to Text-to-Image Diffusion Models.https://arxiv.org/abs/2302.05543 keeps the pretrained backbone and adds trainable control branches for signals like depth, edges, or human pose. This is how you turn "draw a person" into "draw a person in exactly this pose."

Which adaptation tool fits brand lighting, a masked campaign-background repair, and exact pose control?

Answer

Use LoRA for brand-specific lighting or visual style, inpainting for repairing or replacing masked campaign regions, and ControlNet or adapters when the output must follow external structure such as edges, depth, pose, or layout. Don't use generated repairs as factual incident or UI evidence.

Fast inference in practice

DDIM shows the first big shortcut. A production serving pipeline can combine several optimizations, but only after profiling each route. Treat every step count below as a reported or representative route budget, not an SLA. Record checkpoint, resolution, batch or concurrency, hardware, precision, guidance, warmup, end-to-end latency, and quality gates before comparing candidates.

  1. Sampler choice (DDIM / DPM-Solver): DDPM is the canonical stochastic sampler requiring hundreds of steps. DDIM keeps the same training objective but allows deterministic skipped-step trajectories in 20 to 50 steps.[14]Reference 14Denoising Diffusion Implicit Models.https://arxiv.org/abs/2010.02502 Higher-order ODE solvers such as DPM-Solver report high-quality samples in roughly 10 to 20 function evaluations in its experiments. Treat that as paper evidence, not a cross-hardware latency promise.
  2. Latent Consistency Models (LCM): Standard consistency models learn functions that map points on a diffusion ODE trajectory directly to their clean origin x0x_0x0​.[19]Reference 19Consistency Models.https://arxiv.org/abs/2303.01469 Latent Consistency Models (LCM) apply this training in the latent space of pre-trained models such as Stable Diffusion with classifier-free guidance distillation. Their paper reports high-fidelity image generation in 2 to 4 steps.[20]Reference 20Latent Consistency Models: Synthesizing High-Resolution Images with Few-Step Inference.https://arxiv.org/abs/2310.04378 LCM-LoRA provides that acceleration as lightweight plug-in adapters for fine-tuned checkpoints without retraining the base architecture.
  3. Adversarial Diffusion Distillation (ADD / SDXL Turbo): Progressive distillation halves required steps by training student models to match two teacher steps at once.[5]Reference 5Progressive Distillation for Fast Sampling of Diffusion Models.https://arxiv.org/abs/2202.00512 Adversarial Diffusion Distillation (ADD) combines score distillation from a diffusion teacher with an adversarial discriminator loss. Its paper reports single-step to 4-step image synthesis, including SDXL-scale comparisons.[21]Reference 21Adversarial Diffusion Distillation.https://arxiv.org/abs/2311.17042
  4. Straight-line objectives (Rectified Flow / Flow Matching): Rectified flow and flow matching learn transport trajectories between noise and data distributions.[10]Reference 10Flow Straight and Fast: Learning to Generate and Transfer Data with Rectified Flowhttps://arxiv.org/abs/2209.03003[9]Reference 9Flow Matching for Generative Modelinghttps://arxiv.org/abs/2210.02747 Straight paths are intended to make coarse numerical integration effective. Reported 4 to 8-step routes are illustrative; test quality, control, and latency on the target checkpoint and workload.
  5. Guidance and CFG distillation: Distilling classifier-free guidance into a student denoiser can remove separate conditional and unconditional predictions at runtime. If one student prediction replaces two equivalent-quality predictions, denoiser compute changes from 2×2\times2× to 1×1\times1× per step. Verify that equivalence on the release fixture.
  6. Kernel and hardware efficiency: Transformer and attention-heavy denoisers can benefit from fused kernels such as FlashAttention,[22]Reference 22FlashAttention: Fast and Memory-Efficient Exact Attention with IO-Awareness.https://arxiv.org/abs/2205.14135 FP8/INT8 weight quantization, and graph-level optimizations such as TensorRT or torch.compile. Measure memory, warmup, throughput, p50, and p95 because a kernel speedup need not improve end-to-end latency.

In a design-tool setting, those measurements decide which governed routes are affordable. Offline concept boards can tolerate slower samplers; interactive previews may need fewer evaluations and stricter latency budgets. No step count makes output faithful to a real screenshot, so quality and representation checks remain separate gates.

How should a platform team choose between 50-step DDIM and a 6-step distilled model?

Answer

Use the slower sampler when quality, editability, or offline batch generation matters more than latency. Use the distilled model when interactive UX, personalization, or high throughput matters enough to accept quality and control tradeoffs.

Common pitfalls

Use these as symptom-to-owner checks. A visually convincing sample can still fail at guidance, representation, latency, or provenance.

"CFG scale is a probability"

  • Symptom: A team sets w=7.5w=7.5w=7.5 because it "sounds confident," then gets oversaturated textures and repeated details.
  • Cause: www is an extrapolation coefficient between unconditional and conditional noise, not a chance the prompt is followed.
  • Fix: Measure adherence, diversity, and artifacts on the target prompt slice. w=0w=0w=0 is unconditional, w=1w=1w=1 is plain conditional, and larger values overshoot.[13]Reference 13Classifier-Free Diffusion Guidance.https://arxiv.org/abs/2207.12598

"The VAE decoder consumes predicted noise"

  • Symptom: Debugging dumps look like static, or operators think every denoiser step should produce a viewable image.
  • Cause: If the model predicts ϵ\epsilonϵ or vvv, the scheduler converts that into the next latent. Only the final clean latent is decoded.[3]Reference 3Denoising Diffusion Probabilistic Models.https://arxiv.org/abs/2006.11239[6]Reference 6High-Resolution Image Synthesis with Latent Diffusion Models.https://arxiv.org/abs/2112.10752
  • Fix: Inspect ztz_tzt​ after the scheduler update, and decode pixels once at the end.

"Latent diffusion removes the autoencoder bottleneck"

  • Symptom: Fine text, logos, and UI chrome stay mushy even after more denoising steps.
  • Cause: The loop is cheaper because the VAE already discarded high-frequency detail.
  • Fix: Evaluate reconstruction on the details the route must preserve, and don't treat extra sampler steps as a substitute for a better decoder.[6]Reference 6High-Resolution Image Synthesis with Latent Diffusion Models.https://arxiv.org/abs/2112.10752

"DiT is automatically cheaper"

  • Symptom: A transformer backbone is chosen for latency, then token count blows up the attention bill.
  • Cause: DiT scaling results are about quality versus compute, not a promise of cheaper serving.
  • Fix: Patch in latent space, budget tokens, and measure the actual route. Joint text-image attention, as in SD3's MM-DiT, is a published recipe, not a universal speedup.[7]Reference 7Scalable Diffusion Models with Transformers.https://arxiv.org/abs/2212.09748[8]Reference 8Scaling Rectified Flow Transformers for High-Resolution Image Synthesishttps://arxiv.org/abs/2403.03206

"Text diffusion is next-token decoding with a different sampler"

  • Symptom: Serving assumes tokens stream left to right, or treats a missing canvas plan as a tokenizer detail.
  • Cause: Discrete diffusion refines a block. DiffusionGemma is encoder-decoder and block-autoregressive: a 256-token canvas is denoised with bidirectional attention, then appended before the next canvas starts.[1]Reference 1DiffusionGemmahttps://ai.google.dev/gemma/docs/diffusiongemma[12]Reference 12DiffusionGemma Model Cardhttps://ai.google.dev/gemma/docs/diffusiongemma/model_card
  • Fix: Plan canvas length, refinement rounds, commit confidence, and UI behavior for text that can change before it's final.

"A missing watermark proves the image is human-made"

  • Symptom: A copied file without synthetic=true is treated as photographic evidence.
  • Cause: API labels don't travel with the bytes, and detectors can be uncertain after edits.
  • Fix: Embed a watermark, bind the asset digest to the job record, and verify the distributed derivative before publishing. Nondetection still isn't proof of human origin.

Teams building design-asset pipelines can't treat synthetic imagery as verified UI or incident evidence. Generated outputs can show incorrect labels, brand marks, chart values, legal text, or unsafe content.

High-risk routes need output filtering, provenance labeling, watermark verification, and human review before publication. Before running the gate below, predict which outputs should be held: any claim of a real screenshot or use of a logo.

synthetic-output-release-gate.py
1import json 2 3outputs = [ 4 {"asset": "campaign_background", "claims_real_screenshot": False, "contains_logo": False}, 5 {"asset": "ui_mockup_with_text", "claims_real_screenshot": True, "contains_logo": False}, 6 {"asset": "incident_screenshot_evidence", "claims_real_screenshot": True, "contains_logo": True}, 7] 8 9decisions = [] 10for output in outputs: 11 needs_review = output["claims_real_screenshot"] or output["contains_logo"] 12 decisions.append({ 13 "asset": output["asset"], 14 "decision": "human_review" if needs_review else "synthetic_labeled_preview", 15 }) 16 17print(json.dumps(decisions, indent=2))
Output
1[ 2 { 3 "asset": "campaign_background", 4 "decision": "synthetic_labeled_preview" 5 }, 6 { 7 "asset": "ui_mockup_with_text", 8 "decision": "human_review" 9 }, 10 { 11 "asset": "incident_screenshot_evidence", 12 "decision": "human_review" 13 } 14]

Why does synthetic design imagery still need brand-safety checks?

Answer

Generated outputs can misstate UI text, chart values, or reproduce marks and unsafe content. Production systems need output filters, synthetic-content labels, and human review paths for high-risk use cases even when the prompt looks harmless.

Run one generation release audit

A release audit asks whether a candidate improves accepted assets, not whether one sample looks good. Pin model, VAE, scheduler, seed, prompt, negative prompt, step count, and guidance scale.

Vary one axis at a time. Record prompt alignment, text and logo corruption, memorization risk, safety-filter outcomes, latency, and provenance metadata in one release scorecard over the same fixture set.

Use the failure signal to choose the next check. Prompt alignment with garbled text points to representation or VAE limits. Good samples with a p95 breach point to evaluation count, batching, or release work. Uncertain watermark detection blocks publication until the final derivative is checked.

High-impact screenshots, legal text, brand assets, or incident evidence stay behind human review even when automated checks pass.

Complete the lesson

Mastery Check

Answer every question, then check your score. Score 75% or higher to mark this lesson complete.

1.A generated image passes prompt policy and receives a watermark, then a CDN resizes it. The detector checked only the worker's original bytes. When may the job publish its public URL?

Correct answer: After checking the final distributed derivative and binding its digest, model and policy versions, watermark method, and detector result in the durable job record.

The release path must verify the bytes that users will receive, not only an intermediate render. Watermarks are provenance signals rather than universal cryptographic guarantees, and nondetection cannot prove human origin.

2.A trainer uses x0 = 1.0, alpha_bar = 0.5, and epsilon = -0.4 in xt = sqrt(alpha_bar)x0 + sqrt(1-alpha_bar)epsilon. What values and training consequence follow?

Correct answer: xt is about 0.424 and SNR is 1.0, and the closed form samples that timestep directly from x0 and one noise draw.

The signal is about 0.707 and the noise is about -0.283, giving 0.424; SNR is 0.5 divided by 0.5, or 1.0. The marginalized Gaussian transition avoids simulating all earlier forward steps during training.

3.A DDPM training row has true noise [0.20, -0.10, 0.05] and predicted noise [0.18, -0.08, 0.02]. What objective and inference role are correct?

Correct answer: Regress the known added noise with MSE about 0.0006, then let the scheduler use denoiser predictions because x0 is unavailable at sampling.

The squared errors average to about 0.0006, and the standard simple DDPM objective predicts the noise that training injected. At inference, the clean image is unknown, so the learned noise, x0, or velocity prediction guides an approximate reverse update rather than an exact posterior.

4.Using guided = uncond + w × (cond - uncond), uncond = [0.20, -0.10, 0.05], cond = [0.05, -0.35, 0.20], and w = 7.5, which guidance result is correct?

Correct answer: [-0.925, -1.975, 1.175], because the scale extrapolates the conditional-minus-unconditional vector beyond the conditional prediction.

The difference is [-0.15, -0.25, 0.15]. Multiplying by 7.5 and adding the unconditional vector gives [-0.925, -1.975, 1.175]. CFG scale is an extrapolation coefficient, not a probability, and large values can increase artifacts and reduce diversity.

5.A team compares a Stable Diffusion-style U-Net with a DiT backbone for a governed UI illustration route. Which architecture statement is accurate?

Correct answer: A U-Net uses multiscale convolutional down and up paths with skip connections, while a DiT patchifies latent inputs into tokens processed by transformer blocks.

U-Nets preserve local detail through skip paths while building broader context through down and up stages. DiTs change the denoiser backbone to transformer blocks over latent patches. Both still need timestep conditioning, and neither architecture guarantees lower serving cost without measurement.

6.A 512 × 512 × 3 image is generated with a 64 × 64 × 4 continuous latent. Which work reduction and quality caveat should guide the route?

Correct answer: The loop processes 16,384 rather than 786,432 scalar values, a 48x reduction, but VAE reconstruction loss still needs evaluation on important details.

The pixel tensor contains 786,432 values and the latent contains 16,384, so the iterative state is 48 times smaller. Latent diffusion reduces loop cost, but lossy decoding can blur text, logos, or UI details and cannot be fixed merely by adding denoising steps.

7.An offline route uses 50 DDIM steps with ordinary classifier-free guidance. An interactive route uses a distilled six-step model with one prediction per step. What budget and rollout conclusion are sound?

Correct answer: The offline route uses 100 denoiser evaluations and the distilled route uses 6; the faster route still needs separate quality, safety, and editability gates.

Ordinary CFG costs two predictions per step, so 50 × 2 = 100. The stated distilled route costs 6 × 1 = 6. Fewer evaluations can improve latency and throughput, but they do not guarantee teacher-level prompt adherence, controls, or safe governed output.

8.A team evaluates DiffusionGemma-style generation for short code-review summaries. Which serving model matches its dependency structure?

Correct answer: Use a discrete token canvas refined over repeated passes, tuning block length, refinement rounds, confidence-based commits, and UI streaming.

The architecture is block-autoregressive discrete diffusion: an encoder prefills context, a decoder refines a bounded token canvas with bidirectional attention, then the finished canvas is appended before the next block. Tokens can change before commit, so length, stopping, confidence, and presentation are serving decisions.

9.A new sampler and model checkpoint are ready for an image-generation preview route. Which rollout sequence provides meaningful rollback evidence?

Correct answer: Run frozen prompt tests, shadow the candidate, canary 1%, then expand; roll back on safety, p95, memorization, or accepted-quality regressions.

A deterministic fixture and seed help compare versions but do not cover live route behavior. Shadowing and staged canaries expose regressions before broad promotion, while safety, p95 latency, memorization, failure rate, human preference, and cost remain separate release signals.

9 questions remaining.

Next Step
Continue to Real-Time Voice AI Agent

You can now size iterative generation by denoiser evaluations, seeds, and release gates. Next you'll apply that same latency-and-recovery discipline to live audio, where streaming stages, interruptions, and first-audio timing shape the architecture.

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References

DiffusionGemma

Google · 2026

https://ai.google.dev/gemma/docs/diffusiongemma

SynthID

Google DeepMind · 2026

https://deepmind.google/models/synthid/

Denoising Diffusion Probabilistic Models.

Ho, J., Jain, A., & Abbeel, P. · 2020 · NeurIPS 2020

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Progressive Distillation for Fast Sampling of Diffusion Models.

Salimans, T., & Ho, J. · 2022 · ICLR 2022

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High-Resolution Image Synthesis with Latent Diffusion Models.

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https://arxiv.org/abs/2112.10752

Scalable Diffusion Models with Transformers.

Peebles, W., & Chen, S. · 2023 · ICCV 2023

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Scaling Rectified Flow Transformers for High-Resolution Image Synthesis

Esser, P., Kulal, S., Blattmann, A., et al. · 2024

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Flow Matching for Generative Modeling

Lipman, Y., Chen, R. T. Q., Ben-Hamu, H., Nickel, M., & Le, M. · 2022

https://arxiv.org/abs/2210.02747

Flow Straight and Fast: Learning to Generate and Transfer Data with Rectified Flow

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Consistency Models.

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Latent Consistency Models: Synthesizing High-Resolution Images with Few-Step Inference.

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https://arxiv.org/abs/2205.14135

Discussion

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